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Stable Diffusion v1.5

A latent diffusion-based text-to-image generation model that produces photorealistic images from text prompts. It builds upon the Stable Diffusion v1.2 weights and is fine-tuned for improved classifier-free guidance. It can be used via the Diffusers library, ComfyUI, and other interfaces.

Key Information

  • Category: Image Models
  • Source: Huggingface
  • Tags: text-to-image
  • Last updated: January 09, 2026

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Links

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Canonical source: https://huggingface.co/stable-diffusion-v1-5/stable-diffusion-v1-5

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